From the Google AI blog. Interesting ways to improve images based on AI models. Illustrative images at the link.
High Fidelity Image Generation Using Diffusion Models
Friday, July 16, 2021, Posted by Jonathan Ho, Research Scientist and Chitwan Saharia, Software Engineer, Google Research, Brain Team
Natural image synthesis is a broad class of machine learning (ML) tasks with wide-ranging applications that pose a number of design challenges. One example is image super-resolution, in which a model is trained to transform a low resolution image into a detailed high resolution image (e.g., RAISR). Super-resolution has many applications that can range from restoring old family portraits to improving medical imaging systems. Another such image synthesis task is class-conditional image generation, in which a model is trained to generate a sample image from an input class label. The resulting generated sample images can be used to improve performance of downstream models for image classification, segmentation, and more.
Generally, these image synthesis tasks are performed by deep generative models, such as GANs, VAEs, and autoregressive models. Yet each of these generative models has its downsides when trained to synthesize high quality samples on difficult, high resolution datasets. For example, GANs often suffer from unstable training and mode collapse, and autoregressive models typically suffer from slow synthesis speed.
Alternatively, diffusion models, originally proposed in 2015, have seen a recent revival in interest due to their training stability and their promising sample quality results on image and audio generation. Thus, they offer potentially favorable trade-offs compared to other types of deep generative models. Diffusion models work by corrupting the training data by progressively adding Gaussian noise, slowly wiping out details in the data until it becomes pure noise, and then training a neural network to reverse this corruption process. Running this reversed corruption process synthesizes data from pure noise by gradually denoising it until a clean sample is produced. This synthesis procedure can be interpreted as an optimization algorithm that follows the gradient of the data density to produce likely samples.
Today we present two connected approaches that push the boundaries of the image synthesis quality for diffusion models — Super-Resolution via Repeated Refinements (SR3) and a model for class-conditioned synthesis, called Cascaded Diffusion Models (CDM). We show that by scaling up diffusion models and with carefully selected data augmentation techniques, we can outperform existing approaches. Specifically, SR3 attains strong image super-resolution results that surpass GANs in human evaluations. CDM generates high fidelity ImageNet samples that surpass BigGAN-deep and VQ-VAE2 on both FID score and Classification Accuracy Score by a large margin. ...'
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